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--- |
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tags: |
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- text-to-image |
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- stable-diffusion |
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- lora |
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- diffusers |
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- template:sd-lora |
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base_model: stabilityai/stable-diffusion-xl-base-1.0 |
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license: cc-by-nc-nd-4.0 |
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inference: False |
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--- |
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# ⚡ Flash Diffusion: FlashSDXL ⚡ |
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Flash Diffusion is a diffusion distillation method proposed in [Flash Diffusion: Accelerating Any Conditional |
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Diffusion Model for Few Steps Image Generation](http://arxiv.org/abs/2406.02347) *by Clément Chadebec, Onur Tasar, Eyal Benaroche, and Benjamin Aubin* from Jasper Research. |
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This model is a **108M LoRA** distilled version of [SDXL](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) model that is able to generate images in **4 steps**. The main purpose of this model is to reproduce the main results of the paper. |
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See our [live demo](https://huggingface.co/spaces/jasperai/FlashPixart) and official [Github repo](https://github.com/gojasper/flash-diffusion). |
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<p align="center"> |
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<img style="width:700px;" src="images/flash_sdxl.jpg"> |
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</p> |
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# How to use? |
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The model can be used using the `DiffusionPipeline` from `diffusers` library directly. It can allow reducing the number of required sampling steps to **4 steps**. |
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```python |
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from diffusers import DiffusionPipeline, LCMScheduler |
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adapter_id = "jasperai/flash-sdxl" |
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pipe = DiffusionPipeline.from_pretrained( |
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"stabilityai/stable-diffusion-xl-base-1.0", |
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use_safetensors=True, |
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) |
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pipe.scheduler = LCMScheduler.from_pretrained( |
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"stabilityai/stable-diffusion-xl-base-1.0", |
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subfolder="scheduler", |
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timestep_spacing="trailing", |
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) |
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pipe.to("cuda") |
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# Fuse and load LoRA weights |
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pipe.load_lora_weights(adapter_id) |
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pipe.fuse_lora() |
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prompt = "A raccoon reading a book in a lush forest." |
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image = pipe(prompt, num_inference_steps=4, guidance_scale=0).images[0] |
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``` |
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<p align="center"> |
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<img style="width:400px;" src="images/raccoon.png"> |
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</p> |
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# How to use in Comfy? |
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To use FlashSDXL locally using Comfyui you need to : |
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1. Make sure your comfyUI install is up to date |
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2. Download the checkpoint from [huggingface](https://huggingface.co/jasperai/flash-sdxl). |
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In case you wonder how, go to "Files and Version" go to `comfy/` folder and hit the download button next to the `FlashSDXL.safetensors` |
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3. Move the new checkpoint file to your local `comfyUI/models/loras/.` folder |
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4. Use it as a LoRA on top of `sd_xl_base_1.0_0.9vae.safetensors`, a simple comfyui `workflow.json` is provided in this repo (available in the same `comfy/` folder) |
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> Disclaimer : Model has been trained to work with a cfg scale of 1 and a lcm scheduler but parameters can be tweaked a bit. |
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# Combining Flash Diffusion with Existing LoRAs 🎨 |
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FlashSDXL can also be combined with existing LoRAs to unlock few steps generation in a **training free** manner. It can be integrated straight to Hugging Face pipelines. See an example below. |
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```python |
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from diffusers import DiffusionPipeline, LCMScheduler |
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import torch |
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user_lora_id = "TheLastBen/Papercut_SDXL" |
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trigger_word = "papercut" |
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flash_lora_id = "jasperai/flash-sdxl" |
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# Load Pipeline |
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pipe = DiffusionPipeline.from_pretrained( |
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"stabilityai/stable-diffusion-xl-base-1.0", |
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variant="fp16" |
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) |
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# Set scheduler |
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pipe.scheduler = LCMScheduler.from_config( |
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pipe.scheduler.config |
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) |
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# Load LoRAs |
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pipe.load_lora_weights(flash_lora_id, adapter_name="flash") |
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pipe.load_lora_weights(user_lora_id, adapter_name="lora") |
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pipe.set_adapters(["flash", "lora"], adapter_weights=[1.0, 1.0]) |
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pipe.to(device="cuda", dtype=torch.float16) |
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prompt = f"{trigger_word} a cute corgi" |
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image = pipe( |
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prompt, |
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num_inference_steps=4, |
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guidance_scale=0 |
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).images[0] |
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``` |
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<p align="center"> |
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<img style="width:400px;" src="images/corgi.jpg"> |
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</p> |
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> Hint 💡 : You can also use additional LoRA using the provided comfy workflow and test it on your machine. |
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# Combining Flash Diffusion with Existing ControlNets 🎨 |
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FlashSDXL can also be combined with existing ControlNets to unlock few steps generation in a **training free** manner. It can be integrated straight to Hugging Face pipelines. See an example below. |
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```python |
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import torch |
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import cv2 |
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import numpy as np |
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from PIL import Image |
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from diffusers import StableDiffusionXLControlNetPipeline, ControlNetModel, LCMScheduler |
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from diffusers.utils import load_image, make_image_grid |
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flash_lora_id = "jasperai/flash-sdxl" |
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image = load_image( |
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"https://hf.co/datasets/huggingface/documentation-images/resolve/main/diffusers/input_image_vermeer.png" |
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).resize((1024, 1024)) |
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image = np.array(image) |
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image = cv2.Canny(image, 100, 200) |
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image = image[:, :, None].repeat(3, 2) |
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canny_image = Image.fromarray(image) |
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# Load ControlNet |
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controlnet = ControlNetModel.from_pretrained( |
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"diffusers/controlnet-canny-sdxl-1.0", |
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torch_dtype=torch.float16, |
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variant="fp16" |
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) |
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pipe = StableDiffusionXLControlNetPipeline.from_pretrained( |
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"stabilityai/stable-diffusion-xl-base-1.0", |
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controlnet=controlnet, |
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torch_dtype=torch.float16, |
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safety_checker=None, |
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variant="fp16" |
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).to("cuda") |
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# Set scheduler |
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pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config) |
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# Load LoRA |
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pipe.load_lora_weights(flash_lora_id) |
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pipe.fuse_lora() |
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image = pipe( |
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"picture of the mona lisa", |
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image=canny_image, |
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num_inference_steps=4, |
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guidance_scale=0, |
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controlnet_conditioning_scale=0.5, |
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cross_attention_kwargs={"scale": 1}, |
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).images[0] |
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make_image_grid([canny_image, image], rows=1, cols=2) |
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``` |
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<p align="center"> |
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<img style="width:400px;" src="images/controlnet.jpg"> |
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</p> |
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# Training Details |
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The model was trained for 20k iterations on 4 H100 GPUs (representing approximately a total of 176 GPU hours of training). Please refer to the [paper](http://arxiv.org/abs/2406.02347) for further parameters details. |
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**Metrics on COCO 2014 validation (Table 3)** |
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- FID-10k: 21.62 (4 NFE) |
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- CLIP Score: 0.327 (4 NFE) |
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## Citation |
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If you find this work useful or use it in your research, please consider citing us |
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```bibtex |
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@misc{chadebec2024flash, |
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title={Flash Diffusion: Accelerating Any Conditional Diffusion Model for Few Steps Image Generation}, |
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author={Clement Chadebec and Onur Tasar and Eyal Benaroche and Benjamin Aubin}, |
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year={2024}, |
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eprint={2406.02347}, |
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archivePrefix={arXiv}, |
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primaryClass={cs.CV} |
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} |
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``` |
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## License |
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This model is released under the the Creative Commons BY-NC license. |
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