metadata

base_model: stabilityai/stable-diffusion-xl-base-1.0
library_name: diffusers
license: openrail++
tags:
  - text-to-image
  - text-to-image
  - diffusers-training
  - diffusers
  - lora
  - template:sd-lora
  - stable-diffusion-xl
  - stable-diffusion-xl-diffusers
instance_prompt: a photo of sks dog
widget:
  - text: A photo of sks dog in a bucket
    output:
      url: image_0.png
  - text: A photo of sks dog in a bucket
    output:
      url: image_1.png
  - text: A photo of sks dog in a bucket
    output:
      url: image_2.png
  - text: A photo of sks dog in a bucket
    output:
      url: image_3.png

SDXL LoRA DreamBooth - nanaj/flux-lora

Model description

These are nanaj/flux-lora LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0.

The weights were trained using DreamBooth.

LoRA for the text encoder was enabled: False.

Special VAE used for training: madebyollin/sdxl-vae-fp16-fix.

Trigger words

You should use a photo of sks dog to trigger the image generation.

Download model

Weights for this model are available in Safetensors format.

Download them in the Files & versions tab.

Intended uses & limitations

How to use

# TODO: add an example code snippet for running this diffusion pipeline

Limitations and bias

[TODO: provide examples of latent issues and potential remediations]

Training details

[TODO: describe the data used to train the model]