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import gradio as gr |
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from diffusers import StableDiffusionPipeline |
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import torch |
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from huggingface_hub import HfFolder |
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def generate_image(prompt, guidance_scale, num_steps, lora_scale): |
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pipe = StableDiffusionPipeline.from_pretrained( |
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"runwayml/stable-diffusion-v1-5", |
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torch_dtype=torch.float16 |
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).to("cuda") |
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pipe.load_lora_weights( |
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"Scalino84/my-flux-face-v2", |
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weight_name="flux_train_replicate.safetensors" |
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) |
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image = pipe( |
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prompt=prompt, |
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num_inference_steps=num_steps, |
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guidance_scale=guidance_scale, |
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cross_attention_kwargs={"scale": lora_scale} |
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).images[0] |
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return image |
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with gr.Blocks() as demo: |
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gr.Markdown("# Flux Face Generator") |
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with gr.Row(): |
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with gr.Column(): |
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prompt = gr.Textbox( |
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label="Prompt", |
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value="a photo of xyz person, professional headshot", |
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lines=3 |
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) |
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guidance = gr.Slider( |
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label="Guidance Scale", |
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minimum=1, |
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maximum=20, |
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value=7.5, |
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step=0.5 |
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) |
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steps = gr.Slider( |
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label="Inference Steps", |
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minimum=20, |
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maximum=100, |
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value=30, |
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step=1 |
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) |
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lora_scale = gr.Slider( |
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label="LoRA Scale", |
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minimum=0.1, |
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maximum=1.0, |
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value=0.8, |
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step=0.1 |
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) |
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generate = gr.Button("Generate Image") |
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with gr.Column(): |
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output = gr.Image(label="Generated Image") |
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generate.click( |
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fn=generate_image, |
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inputs=[prompt, guidance, steps, lora_scale], |
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outputs=output |
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) |
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demo.launch() |
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